Stable diffusion face refiner online reddit

0 and upscalers. The title tells everything. It works to a degree but maybe not enough. Key Takeaways. It depends on: what model you are using for the refiner (hint, you don't HAVE to use stabilities refiner model, you can use any model that is the same family as the base generation model - so for example a SD1. I was surprised by how nicely the SDXL Refiner can work even with Dreamshaper as long as you keep the steps really low. I created this comfyUI workflow to use the new SDXL Refiner with old models: json here. 3), detailed face, /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app So in order to get some answers I'm comparing SDXL1. Hires fix is still there, you just need to click to expand but face restore has indeed been removed from the main page. This simple thing made me a fan of Stable Diffusion. Whenever you generate images that have a lot of detail and different topics in them, SD struggles to not mix those details into every "space" it's filling in running through the denoising step. 0 includes the following experimental functions: Free Lunch (v1 and v2) AI researchers have discovered an optimization for Stable Diffusion models that improves the quality of the generated images. I will see that when I click on the wrong model and used it instead of the base. Edit: RTX 3080 10gb example with a shitty prompt just for demonstration purposes: Without --medvram-sdxl enabled, base SDXL + refiner took 5 mins 6. F222 is a traditional fine tuned model that does not require a special keyword. Try: rear view shot or just rear shot. ComfyUI with SDXL (Base+Refiner) + ControlNet XL OpenPose + FaceDefiner (2x) ComfyUI is hard. 3) Jul 22, 2023 · After Detailer (adetailer) is a Stable Diffusion Automatic11111 web-UI extension that automates inpainting and more. I've heard you get better results with full body shots if the source images used for the training were also full body shots, and also keeping the dimension to no more than 512X512 durign generation. Unrealistic body standards having a 24 pack. 9 (just search in youtube sdxl 0. Hi. Award. Here is an example of two images. I see a lot of people complaining about the new hires. As you can see the difference is an improvement but the image retains nearly everything. [Cross-Post] I feel the original one is better, high denoise refiner destroys the lighting consistency, especially the hands become flat and even changed the skin color in the third image. face recognition API. i tried "camera from behind" or "camera shot from behind", i cant really think of other prompts to use, but iv only been able to get like 1 out of 20 images to be from behind with this. Normal Hires. 0 for ComfyUI - Now with Face Swapper, Prompt Enricher (via OpenAI), Image2Image (single images and batches), FreeU v2, XY Plot, ControlNet and ControlLoRAs, SDXL Base + Refiner, Hand Detailer, Face Detailer, Upscalers, ReVision, etc. For prompt use something like face, eye color, hair color, hair style, expression. Switch the timing point from 0. Im using automatic1111 and I run the initial prompt with sdxl but the lora I made with sd1. For example, if you wanted a great image of a person in a firefighter outfit, you could add a specific extra ‘embedding’ model trained on images of firefighter outfits. Use 0. For today's tutorial I will be using Stable Diffusion XL (SDXL) with the 0. These comparisons are useless without knowing your workflow. It is suitably sized to become the next standard in text-to-image models. Below 0. 9 vae in there. What they have is a marketable product. Reply. 0 Base vs Base+refiner comparison using different Samplers. The refiner was trained in tandem with the base, so it will not work without it. 5 and 2. In summary, it's crucial to make valid comparisons when evaluating the SDXL with and without the refiner. 5 refiner node. On the other hand tin my experience the SD3 renders doesn't mix very well with refiners so what you obtain is almost a dead end. choose two different styles of models, one as a base, one as a refinement of the model. 5 for final detail refinement seems to give me the ultimate control. Stable Diffusion looks too complicated”. AP Workflow 5. Second, you'll need the photo style SDXL checkpoint of your preference. I've found very good results doing 15-20 steps with SDXL which produces a somewhat rough image, then 20 steps at 0. The first is PixArt Sigma with no refinement and the second is after a . This series of images is made to see if the color depth in SD3 can be translated in the process of refiner pass. This brings back memories of the first time that I use Stable Diffusion myself. 0 and some of the current available custom models on civitai with and without the refiner. Stable Diffusion creates images out of pure noise. 7 in the Denoise for Best results. First get a photo of the head in the same direction as the result in mind. I was really like stable Cascade mixed with a 1. Code for automatically detecting and correcting hands in Stable Diffusion using models of hands, ControlNet, and inpainting. People using utilities like Textual Inversion and DreamBooth have been able to solve the problem in narrow use cases, but to the best of my knowledge there isn't yet a reliable solution to make on-model characters without just straight up hand-holding the AI. The control Net Softedge is used to preserve the elements and shape, you can also use Lineart) 3) Setup Animate Diff Refiner Hey, a lot of thanks for this! I had a pretty good face upscaling routine going for 1. realvisXL is great and currently probably better than Juggernaut however it is not a Pony model so it can't do what Pony can (but can do a few things that Pony struggles with, like working with controlnet). You'll need two checkpoints "Real Pony" checkpoint (there are three, you want the one with the most upvotes on civit to start): this will be the primary checkpoint. These sample images were created locally using Automatic1111's web ui, but you can also achieve similar results by entering prompts one at a time into your distribution/website of choice. The negative prompt sounds like frustration to me. It's called Family pack, get with the times old man! Dude is corn. 519K subscribers in the StableDiffusion Make sure you have: Settings -> Stable Diffusion -- > "Maximum number of checkpoints loaded at the same time" set to 2 so it wont unload and reload the model for each pass. 9 refiner node. The vae that was originally baked into SDXL created visual artifacts when it tried to do its "invisible watermarking". (as mentioned Used Automatic1111, SDXL 1. If you have powerful GPU and 32GB of RAM, plenty of disc space - install ComfyUI - snag the workflow - just an image that looks like this one that was made with Comfy - drop it in the UI - and write your prompt - but the setup is a bit involved - and things don't always go smoothly - you will need the toon model as well - Civitai/HuggingFace I can't get Outpainting to work in Stable Diffusion. The model doesn’t seem to work for anime images…. They also have an SDXL Lora that kinda adds some contrast. It is the delicate interplay of shadow and light. 5 VAE) -> SD 1. There is no such thing as an SD 1. It'll be at the top though, not where it used to be. To get it back, go to settings --> user interface and add it back. Misconfiguring nodes can lead to erroneous conclusions, and it's essential to understand the correct settings for a fair assessment. Forcing Lora weights higher breaks the ability for generalising pose, costume, colors, settings etc. 01 ~ 1, each increase of 0. 0 Base, moved it to img2img, removed the LORA and changed the checkpoint to SDXL 1. But they have different CFG values (because of the lightning). img2img API with inpainting. If you aren't using that GUI, the best option is to bring it into GIMP Fooocus-MRE v2. Uncharacteristically, it's not as tidy as I'd like, mainly due to a challenge I have with passing the checkpoint/model name through reroute nodes. So, I'm mostly getting really good results in automatic1111 Yes its human faces only, probably best prompting on your dogs photo using img2img or controlnet. Do not use the high res fix section (can select none, 0 steps in the high res section), go to the refiner section instead that will be new with all your other extensions (like control net or whatever other extensions you Using refiner with different settings. I haven't played with Dreambooth myself so just going by other people's experience. Can take a while, on average I need two or the dalle2 inpainting prompts to get them fixed. For now, I have to manually copy the right prompts. I'll then be wondering why the image was so bad 😂. and have to close terminal and restart a1111 again to clear that OOM effect. Start with a denoise around . 5 version, losing most of the XL elements. 70 Prompt Comparison: SD3 API vs SD3 Medium. At that moment, I was able to just download a zip, type something in webui, and then click generate. First of all, sorry if this doesn't make sense, i'm french so english isn't my native language and i'm self-taught when it comes to english. But I'm not sure what I'm doing wrong, in the controlnet area I find the hand depth model and can use it, I would also like to use it in the adetailer (as described in Git) but can't find or select the depth model (control_v11f1p_sd15_depth) there. Same with SDXL, you can use any two SDXL models as the base Hopefully Adetailer gets updated soon so you can choose the hands inpaint model instead of inpaint global harmonious. In my opinion the renders of pixart tend to be more interesting and beautiful than the SD3 renders, but they need a second pass with a refiner. 0. Consistent character faces, designs, outfits, and the like are very difficult for Stable Diffusion, and those are open problems. If you're using Automatic's GUI there should be an option for full res inpainting so you can mask off the face and generate a new one using a prompt referencing the face and it will generate the face at the full resolution of the image and then scale it down to fit the mask. I also automated the split of the diffusion steps between the Base and the I thought my gaming would be at least a lot better than my 2070 super 8gb. Trained information is represented in an alternate form (using CLIP for text and VAE for image embeddings). After some testing I think the degradation is more noticeable with concepts than styles. 3~0. It will just produce distorted, incoherent images. It is the curve of rolling hills. 4 - 0. 1. I will first try out the newest sd. 6. I have tried uninstalling stable diffusion (deleting Taking a good image with a poor face, then cropping into the face at an enlarged resolution of it's own, generating a new face with more detail then using an image editor to layer the new face on the old photo and using img2img again to combine them is a very common and powerful practice. Comparison. 9 safetesnors file. 0 of my AP Workflow for ComfyUI. If the problem still persists I will do the refiner-retraining. And the SDE++ 2M versions are also fast per step. It'll be perfect if it includes upscale too (though I can upscale it in an extra step in the extras tap of automatic1111). This one feels like it starts to have problems before the effect can Is there an explanation for how to use the refiner in ComfyUI? You can just use someone elses workflow of 0. This simple thing also made my that friend a fan of Stable Diffusion. The main reason why I chose to do this is a selfish one. 0 base model and HiresFix x2. In any case, we could compare the picture obtained with the correct workflow and the refiner. 5 denoise with SD1. 5 model in highresfix with denoise set in the . Ah also, death to the false emperor, blood for the blood god. Can someone guide me to the best all-in-one workflow that includes base model, refiner model, hi-res fix, and one LORA. Interesting, gonna try this tomorrow. 📷 All of the flexibility of Stable Diffusion: SDXL is primed for complex image design workflows that include generation for text or base image, inpainting . Yep! I've tried and refiner degrades (or changes) the results. But if you use both together it will make very little differences. Use img2img to refine details. The smaller size of this model makes it perfect for running on consumer PCs and laptops as well as enterprise-tier GPUs. Thanks. 78. In this post, you will learn how it works, how to use it, and some common use cases. No Automatic1111 or ComfyUI node as of yet. , that is more conspicuous than the number of fingers RTX 3060 12GB VRAM, and 32GB system RAM here. 6 or too many steps and it becomes a more fully SD1. 5 can be seen in the style of the obvious changes in 0. No surprises, Medium is much worse. Technical details regarding Stable Diffusion samplers, confirmed by Katherine: - DDIM and PLMS are originally the Latent Diffusion repo DDIM was implemented by CompVis group and was default (slightly different update rule than the samplers below, eqn 15 in DDIM paper is the update rule vs solving eqn 14's ODE directly) /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. I just released version 4. i dont understand what you need. I want to use Pony as a base model and Juggernaut Lightning as a refiner for more realistic images. 9 and Stable Diffusion 1. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Opening the image in stable-diffusion-webui's PNG-info I can see that there are indeed two different sets of prompts in that file and for some reason the wrong one is being chosen. 5 model as your base model, and a second SD1. I've tried this article, but the result does not give me what I wanted. You just can't change the conditioning mask strength like you can with a proper inpainting model, but most people don't even know what that is. 2), cottage. Edit: I realized that the workflow loads just fine, but the prompts are sometimes not as expected. 1 sdxl model and 1 sd1. 1), crowded, alluring eyes, detailed skin, highly detailed, hyperdetailed, intricate, soft lighting, deep focus, photographed on a Canon 5D, 24mm macro lens, F/8 aperture, film still [after]{zoom_enhance mask="face" replacement /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Kohya Deepshrink is based on Scalecrafter research. SD 1. They mostly use python to train Stable Diffusion. All dreambooth models require a special keyword to condition the image generation but the traditional fine tuning (continue training with a narrow dataset) doesn't. It adds detail and cleans up artifacts. Then install the SDXL Demo extension . 74 votes, 16 comments. I had to use clip interrogator on Replicate because it gives me errors when using it locally. This seemed to add more detail all the way up to 0. Medium is using the base workflow on the huggingface Using Refiner -> Base or just CrystalClearXL or other model from the start -> VAEDecode->VAEEncode (SD 1. 9 vae, along with the refiner model. 5 model. When SD tries to generate an image that is too different in size and aspect ratio from what it is trained on, you end up getting elongated or multiple features such as two heads, two torsos, and 4 legs. In my experience, LMS has similar quality to 2M (both Karras and non, compared to their 2M versions), but LMS samplers are more artifact-prone, esp. I just started learning about Stable Diffusion recently, I downloaded the safe-tensors directly from huggingface for Base and Refiner model, I found…. AP Workflow v5. example of the basic model of photo-realistic style, the model used in the refiner is the anime style. next version as it should have the newest diffusers and should be lora compatible for the first time. Basically a bunch of junk so that I can perfect the image. If you use ComfyUI you can instead use the Ksampler Regarding the "switching" there's a problem right now with the 1. 0 model. Thanks tons! Accidentally used the refiner model to generate images. Then I do multiple img2img passes with a higher resolution, more VATSIM (Virtual Air Traffic Simulation Network) is the go-to online flight simulation network, where virtual pilots can connect their flight simulators to a shared network and enjoy realistic communication and procedures by VATSIM's trained virtual Air Traffic Controllers. Steps: (some of the settings I used you can see in the slides) Generate first pass with txt2img The checkboxes (face fix, hires fix) disappeared. As per the SD super stage event, the refiner is an optional second pass that can improve some generations. LewdGarlic. Can I use a different CFG value for the refiner in Comfy? I'm currently using Forge, should I switch to Comfy or StableSwarm? 2. It's often not required. Looking for a tutorial to train your own face using Automatic1111. 5 model as the "refiner"). It detects hands greater than 60x60 pixels in a 512x512 image, fits a mesh model and then generates SDXL vs SDXL Refiner - Img2Img Denoising Plot. Stable Diffusion 3 Medium is Stability AI’s most advanced text-to-image open model yet, comprising two billion parameters. Structured Stable Diffusion courses. This accuracy allows much more to be done to get the perfect image directly from text, even before using the more advanced features or fine-tuning that Stable Diffusion is famous for. 55 and go from there. text_l & refiner: "(pale skin:1. E. Reply reply [deleted] High detail RAW color Photo of a strong man, hands in the face, urban city in the background, (full body view:1. SDXL vs DreamshaperXL Alpha, +/- Refiner. An style can be slightly changed in the refining step, but a concept that doesn't exist in the standard dataset is usually lost or turned into another thing (I. Thanks for this - newbs coming from A1111 can be overwhelmed by the ComfyUI when trying to locate nodes. The refiner should definitely NOT be used as the starting point model for text2img. Generate your images through automatic1111 as always and then, go to the SDXL Demo extension tab, turn on 'Refine' checkbox and drag your image onto the square. 509K subscribers in the StableDiffusion ComfyUI with SDXL (Base+Refiner) + ControlNet XL OpenPose + FaceDefiner (2x) ComfyUI is hard. Here are the solutions: ***Basically, install the refiner extension (sd-webui-refiner). that extension really helps. Refiner extension not doing anything. a close up of a woman with a butterfly on her head, a photorealistic painting, by Anna Dittmann Flexibility. Automatic1111 can’t use the refiner correctly. Setup. Basically it just creates a 512x512 as usual, then upscales it, then feeds it to the refiner. We all know SD web UI and ComfyUI - those are great tools for people who want to make a deep dive into details, customize workflows, use advanced extensions, and so on. 5 LCM refiner sampler pass. 0 Refine. Even the Comfy workflows aren’t necessarily ideal, but they’re at least closer. Yes. We need laws that mark images like this as AI generated so we don't get low self-esteem. fix while using the refiner you will see a huge difference. We'll see about the actual quality, flexibility, prompt adherence and optimization, if/when SD3 comes out fully. So the website shows all the images SD was trained on and more. 40 denoise using Zavy 's excellent ZavyChromaXL v7. So they put some stuff with the . DreamshaperXL is really new so this is just for fun. Create a Load Checkpoint node, in that node select the sd_xl_refiner_0. Karras in general are superb at low step counts (though LMS Karras gets lots of artifacts at high step counts, so never do that). Tried a bunch of images, and none of them had the hands detected…. I know there is the ComfyAnonymous workflow but it's lacking. The soft inpainting feature is also handy, it tends to blend the seams very well on the inpainted area. 5 models and LoRA are so fine tuned that while SDXL gives me a much wider range of control, getting the 'perfect' finish seems to only be reliable with 2) Set Refiner Upscale Value and Denoise value. I did run the prompts on the SD3 API again to make sure they haven't changed it, and the results were the same, so it's still good. but if I run Base model (creating some images with it) without activating that extension or simply forgot to select the Refiner model, and LATER activating it, it gets OOM (out of memory) very much likely when generating images. Imagine if you can do your model photoshoot with your new watch, skin care product, or line of overprice handbags in a studio, and seamlessly put the model in the streets of Milan, on the beaches of the Maldives, or wherever else instagram and tiktok says your target demo wants UNet does precisely this, working on different levels of detail as it downscales and upscales. All prompts share the same seed. 2) and used the following negative - Negative prompt: blurry, low quality, worst quality, low resolution, artifacts, oversaturated, text, watermark, logo, signature, out SDXL 1. 9 workflow, the one that olivio sarikas video works just fine) just replace the models with 1. A TON of budget of commercial shoots is location-based. The prompts: (simple background:1. From L to R, this is SDXL Base -- SDXL + Refiner -- Dreamshaper -- Dreamshaper + SDXL Refiner. I am going to experiment a bit more but if it doesn't work out, I may just use Pixart for the global compositional coherence latent base for SDXL and SD 1. Here's how to get the benefits of Pony XL, without the drawbacks of art-style. You can inpaint with SDXL like you can with any model. I'd like to share Fooocus-MRE (MoonRide Edition), my variant of the original Fooocus (developed by lllyasviel ), new UI for SDXL models. 5 of my wifes face works much better than the ones Ive made with sdxl so I enabled independent prompting(for highresfix and refiner) and use the 1. Inpainting is almost always needed to fix the face consistency. 8. I fix all my hands in dalle-2. Anyway, while i was writing this post, there has been a new update and it now look like this : Here we go. i came across the "Refiner extension" in the comments here described as "the correct way to use refiner with SDXL" but i am getting the exact same image between checking it on and off and generating the same image seed a few times as a test. I have been using automatic1111, don't know much about comfyui. A person face changes after ADMIN MOD. I had the same idea of retraining it with the refiner model and then load the lora for the refiner model with the refiner-trained-lora. I also automated the split of the diffusion steps between the Base and the Actually the normal XL BASE model is better than the refiner in some points (face for instance) but I think that the refiner can bring some interesting details. But stable diffusion is faster and I can load the workflow I like. Like there's Embeddings, which there are quite a few Me too! realvisXL is awesome at photorealism. 5 in A1111, but now with SDXL in Comfy I'm struggling to get good results by simply sending an upscaled output to a new pair of base+refiner samplers Code Posted for Hand Refiner. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. It's not, it's just barely better. 7 in the Refiner Upscale to give a little room in the image to add details. Install the SDXL auto1111 branch and get both models from stability ai (base and refiner). Note: I used a 4x upscaling model which produces a 2048x2048, using a 2x model should get better times, probably with the same effect. 0. I guess an important thing for the quality (when you mention without finetunes) is that this time, the base model is finetuned by Lykon, the number 1 model creator on civitai. Use a value around 1. Color problems when using face reference. 2), (isometric 3d art of floating rock citadel:1), cobblestone, flowers, verdant, stone, moss, fish pool, (waterfall:1. Not sure if it’s the quality of the image or something but the colors become horrible and the art style becomes much more stylized instead of realistic, even when I try higher resolution images. fix @Dr__Macabre. Is there currently a way to adapt a ComfyUI workflow to avoid the refiner touching any human faces? It's removing details that I want kept there: it makes all faces smooth, de-ages them (I don't want that!) and evens them out, which deletes all of the characters' personalities, age, and uniqueness as a result. Question - Help. But it's reasonably clean to be used as a learning tool, which is and will It's amazing - I can get 1024x1024 SDXL images in ~40 seconds at 40 iterations euler A with base/refiner with the medvram-sdxl flag enabled now. I've search but found nothing that seems to use Automatic1111. 45 denoise it fails to actually refine it. Simply ran the prompt in txt2img with SDXL 1. It is not a reasonable approximation, it is the actual data it was trained on. fix with SDXL is broken. Last, I also performed the same test with a resize by scale of 2: SDXL vs SDXL Refiner - 2x Img2Img Denoising Plot. The dataset I linked above contains 5 billion images, it's called LAION-5B. Suppose we want a bar-scene from dungeons and dragons, we might prompt for something like. 0, VAE hash example of workflow: Prompt: full body photo of beautiful age 18 girl, elf ears, blonde hair, freckles, sexy, beautiful, BREAK hiding behind a tree in the forest /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. SD is a big thing with a lot going on, don't be afraid The truth about hires. Ideally the refiner should be applied at the generation phase, not the upscaling phase. 5, we're starting small and I'll take you along the entire journey. This is not my code, I'm simply posting it. Just like Juggernaut started with Stable Diffusion 1. Legal and PR issues. Every time I use a face reference for my stable diffusion model I get really weird artifacts. If you don't use hires. Not too sure how exactly to do all this others that are up to date will know better. 0-RC. 7 and then close to the base model. Describe the character and add to the end of the prompt: illustration by (Studio ghibli style, Art by Hayao Miyazaki:1. Workflow Overview: txt2Img API. Use SDXL Refiner with old models. 2. 30ish range and it fits her face lora to the image without Very nice. Experimental Functions. If you install locally you can add in your own additions to the main model. 5. 5 secs More than 0. If you've seen this post before, you know what to expect. I downscale my SD pictures before using them in dalle-2, then do img2img again and work with cfg and init strength till they just retouch the dalle2 hands. It can do this because it was trained on a lot of images paired with their text captioning with various amount of noise added to the image. It saves you time and is great for quickly fixing common issues like garbled faces. Stable Diffusion is trained on a subset of those images, around 600 million of those, supposedly. However, this also means that the beginning might be a bit rough ;) NSFW (Nude for example) is possible, but it's not yet recommended and can be prone to errors. 85, although producing some weird paws on some of the steps. After an entire weekend reviewing the material, I think (I hope!) I got the implementation right: As the title says, I included ControlNet XL OpenPose and FaceDefiner models. I'm using the recommended settings; Sampling Steps: 80-100, Sampler: Euler a, Denoising strength: 0. (viewed from behind:1. 3 - 1. Then play with the refiner steps and strength (30/50 Activate the Face Swapper via the auxiliary switch in the Functions section of the workflow. There is an SDXL 0. Step two - upscale: Change the model from the SDXL base to the refiner and process the raw picture in img2img using the Ultimate SD upscale extension with the following settings: (same prompt) Steps: 30, Sampler: DPM++ 2M Karras, CFG scale: 7, Seed: 1799556987, Size: 2304x1792, Model hash: 7440042bbd, Model: sd_xl_refiner_1. pony is anime and 2d model. fix. Technically dreambooth is a also a fine tuning technique. pm ve tz xl xz jy rj bq kh by